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FLUX1.1 [pro] ultra
FLUX1.1 [pro] ultra is the newest version of FLUX1.1 [pro], maintaining professional-grade image quality while delivering up to 2K resolution with improved photo realism.
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FLUX1.1 [pro] ultra is the newest version of FLUX1.1 [pro], maintaining professional-grade image quality while delivering up to 2K resolution with improved photo realism.
Train styles, people and other subjects at blazing speeds.
Recraft V3 is a text-to-image model with the ability to generate long texts, vector art, images in brand style, and much more. As of today, it is SOTA in image generation, proven by Hugging Face’s industry-leading Text-to-Image Benchmark by Artificial Analysis.
Generate video clips from your images using MiniMax Video model
AuraFlow v0.3 is an open-source flow-based text-to-image generation model that achieves state-of-the-art results on GenEval. The model is currently in beta.
FLUX.1 Image-to-Image is a high-performance endpoint for the FLUX.1 [dev] model that enables rapid transformation of existing images, delivering high-quality style transfers and image modifications with the core FLUX capabilities.
Super fast endpoint for the FLUX.1 [dev] model with LoRA support, enabling rapid and high-quality image generation using pre-trained LoRA adaptations for personalization, specific styles, brand identities, and product-specific outputs.
OmniGen is a unified image generation model that can generate a wide range of images from multi-modal prompts. It can be used for various tasks such as Image Editing, Personalized Image Generation, Virtual Try-On, Multi Person Generation and more!
Stable Diffusion 3.5 Large is a Multimodal Diffusion Transformer (MMDiT) text-to-image model that features improved performance in image quality, typography, complex prompt understanding, and resource-efficiency.
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Explore all available models provided by fal.ai
FLUX.1 [dev] is a 12 billion parameter flow transformer that generates high-quality images from text. It is suitable for personal and commercial use.
FLUX.1 [schnell] is a 12 billion parameter flow transformer that generates high-quality images from text in 1 to 4 steps, suitable for personal and commercial use.
FLUX1.1 [pro] is an enhanced version of FLUX.1 [pro], improved image generation capabilities, delivering superior composition, detail, and artistic fidelity compared to its predecessor.
FLUX.1 [pro] new is an accelerated version of FLUX.1 [pro], maintaining professional-grade image quality while delivering significantly faster generation speeds.
Stable Diffusion 3.5 Medium is a Multimodal Diffusion Transformer (MMDiT) text-to-image model that features improved performance in image quality, typography, complex prompt understanding, and resource-efficiency.
Recraft V3 Create Style is capable of creating unique styles for Recraft V3 based on your images.
FLUX Realism LoRA is a specialized fine-tuning adaptation that enhances FLUX models to produce hyper-realistic images with exceptional detail, accurate lighting, and true-to-life textures. Optimized for photographic quality and real-world accuracy.
FLUX LoRA Inpainting is a specialized endpoint that enables precise image editing and completion with FLUX models, combining rapid processing speeds with LoRA adaptation capabilities for high-quality selective image generation and modification.
FLUX LoRA Image-to-Image is a high-performance endpoint that transforms existing images using FLUX models, leveraging LoRA adaptations to enable rapid and precise image style transfer, modifications, and artistic variations.
A versatile endpoint for the FLUX.1 [dev] model that supports multiple AI extensions including LoRA, ControlNet conditioning, and IP-Adapter integration, enabling comprehensive control over image generation through various guidance methods.
FLUX General Inpainting is a versatile endpoint that enables precise image editing and completion, supporting multiple AI extensions including LoRA, ControlNet, and IP-Adapter for enhanced control over inpainting results and sophisticated image modifications.
FLUX General Image-to-Image is a versatile endpoint that transforms existing images with support for LoRA, ControlNet, and IP-Adapter extensions, enabling precise control over style transfer, modifications, and artistic variations through multiple guidance methods.
A specialized FLUX endpoint combining differential diffusion control with LoRA, ControlNet, and IP-Adapter support, enabling precise, region-specific image transformations through customizable change maps.
A general purpose endpoint for the FLUX.1 [dev] model, implementing the RF-Inversion pipeline. This can be used to edit a reference image based on a prompt.
An endpoint for personalized image generation using Flux as per given description.
An endpoint for re-lighting photos and changing their backgrounds per a given description
FLUX.1 Differential Diffusion is a rapid endpoint that enables swift, granular control over image transformations through change maps, delivering fast and precise region-specific modifications while maintaining FLUX.1 [dev]'s high-quality output.
Stable Diffusion 3 Medium (Text to Image) is a Multimodal Diffusion Transformer (MMDiT) model that improves image quality, typography, prompt understanding, and efficiency.
Stable Diffusion 3 Medium (Image to Image) is a Multimodal Diffusion Transformer (MMDiT) model that improves image quality, typography, prompt understanding, and efficiency.
Run SDXL at the speed of light
Run Any Stable Diffusion model with customizable LoRA weights.
Upscale your images with AuraSR.
Stable Cascade: Image generation on a smaller & cheaper latent space.
Generate video clips from your prompts using MiniMax model
Transform text into hyper-realistic videos with Haiper 2.0. Experience industry-leading resolution, fluid motion, and rapid generation for stunning AI videos.
Transform text into hyper-realistic videos with Haiper 2.0. Experience industry-leading resolution, fluid motion, and rapid generation for stunning AI videos.
Mochi 1 preview is an open state-of-the-art video generation model with high-fidelity motion and strong prompt adherence in preliminary evaluation.
Generate video clips from your prompts using Luma Dream Machine v1.5
Generate video clips from your images using Luma Dream Machine v1.5
Generate video clips from your prompts using Kling 1.0
Generate video clips from your images using Kling 1.0
Generate video clips from your prompts using Kling 1.0 (pro)
Generate video clips from your images using Kling 1.0 (pro)
Generate videos from prompts using CogVideoX-5B
Generate videos from videos and prompts using CogVideoX-5B
Generate videos from images and prompts using CogVideoX-5B
Generate short video clips from your images using SVD v1.1
Generate short video clips from your prompts using SVD v1.1
Generate short video clips from your images using SVD v1.1 at Lightning Speed
bilateral reference framework (BiRefNet) for high-resolution dichotomous image segmentation (DIS)
Generate short video clips from your images using SVD v1.1 at Lightning Speed
Create creative upscaled images.
Clarity upscaler for images with high fidelity.
SOTA Image Upscaler
Run SDXL at the speed of light
Run SDXL at the speed of light
Run SDXL at the speed of light
Run SDXL at the speed of light
Run SDXL at the speed of light
Run SDXL at the speed of light
Whisper is a model for speech transcription and translation.
[Experimental] Whisper v3 Large -- but optimized by our inference wizards. Same WER, double the performance!
Run SDXL at the speed of light
Run SDXL at the speed of light
Run SDXL at the speed of light
Hyper-charge SDXL's performance and creativity.
Hyper-charge SDXL's performance and creativity.
Hyper-charge SDXL's performance and creativity.
State-of-the-art open-source model in aesthetic quality
State-of-the-art open-source model in aesthetic quality
State-of-the-art open-source model in aesthetic quality
Interpolate between video frames
Interpolate between image frames
Generate short video clips from your prompts
SD 1.5 ControlNet
Customizing Realistic Human Photos via Stacked ID Embedding
Produce high-quality images with minimal inference steps.
Produce high-quality images with minimal inference steps. Optimized for 512x512 input image size.
Default parameters with automated optimizations and quality improvements.
Re-animate your videos with evolved consistency!
Re-animate your videos with evolved consistency!
Animate your ideas!
Re-animate your videos!
Animate your ideas in lightning speed!
Re-animate your videos in lightning speed!
Create illusions conditioned on image.
Create depth maps using Midas depth estimation.
Remove the background from an image.
Upscale images by a given factor.
Generate Images with ControlNet.
Generate Images with ControlNet.
Generate Images with ControlNet.
Inpaint images with SD and SDXL
Animate Your Drawings with Latent Consistency Models!
Tuning-free ID customization.
High quality zero-shot personalization
Create depth maps using Marigold depth estimation.
Open source text-to-audio model.
Diffusion based high quality edge detection
State of the art Image to 3D Object generation
Default parameters with automated optimizations and quality improvements.
Default parameters with automated optimizations and quality improvements.
Default parameters with automated optimizations and quality improvements.
Automatically retouches faces to smooth skin and remove blemishes.
Use any large language model from our selected catalogue (powered by OpenRouter)
Use any vision language model from our selected catalogue (powered by OpenRouter)
Vision
Vision
Predict the probability of an image being NSFW.
Fooocus extreme speed mode as a standalone app.
Fooocus extreme speed mode as a standalone app.
Create stickers from faces.
Answer questions from the images.
Learning Realistic 3D Motion Coefficients for Stylized Audio-Driven Single Image Talking Face Animation
MuseTalk is a real-time high quality audio-driven lip-syncing model. Use MuseTalk to animate a face with your own audio.
Learning Realistic 3D Motion Coefficients for Stylized Audio-Driven Single Image Talking Face Animation
SDXL with an alpha channel.
Stable Diffusion v1.5
Run Any Stable Diffusion model with customizable LoRA weights.
Run SDXL at the speed of light
Run SDXL at the speed of light
Run Any Stable Diffusion model with customizable LoRA weights.
Weak-to-Strong Training of Diffusion Transformer for 4K Text-to-Image Generation
Dreamshaper model.
Generate realistic images.
Collection of SDXL Lightning models.
Any pose, any style, any identity
Image based Virtual Try-On
Predict poses.
Anime finetune of Würstchen V3.
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
Florence-2 is an advanced vision foundation model that uses a prompt-based approach to handle a wide range of vision and vision-language tasks
A powerful image to novel multiview model with normals.
Transfer expression from a video to a portrait.
Transfer expression from a video to a portrait.
Photorealistic Text-to-Image
An efficent SDXL multi-controlnet text-to-image model.
An efficent SDXL multi-controlnet image-to-image model.
An efficent SDXL multi-controlnet inpainting model.
SAM 2 is a model for segmenting images and videos in real-time.
SAM 2 is a model for segmenting images and videos in real-time.
Multimodal vision-language model for single/multi image understanding
Multimodal vision-language model for video understanding
Animate a reference image with a driving video using ControlNeXt.
Various image preprocessing tools for ControlNet and other applications.
Canny edge detection preprocessor.
Depth Anything v2 preprocessor.
Holistically-Nested Edge Detection (HED) preprocessor.
Line art preprocessor.
MiDaS depth estimation preprocessor.
M-LSD line segment detection preprocessor.
PIDI (Pidinet) preprocessor.
Segment Anything Model (SAM) preprocessor.
Scribble preprocessor.
TEED (Temporal Edge Enhancement Detection) preprocessor.
ZoeDepth preprocessor.
F5 TTS